40 stories tagged diffusion-models, newest first.
27 points by matt_d | 2 comments
Research blog post introducing OPAL, a reactive optimizer where an LLM agent monitors training dynamics and rewrites the optimizer's source code during a training run.
270 points by westsmith | 88 comments
A personal guide to AI image generation covering history, how diffusion and multimodal models work, costs, prompting, and integrating image APIs into code and agents.
201 points by gmays | 57 comments
Black Forest Labs blog post announcing FLUX 3 Image Klein, a 5.6B-parameter open-weights image generation model.
330 points by dvrp | 106 comments
Product launch page for Krea 3, Krea's frontier AI image editing model supporting region-based prompted edits at scale.
173 points by jsomers | 53 comments
Jane Street blog post exploring whether autoregressive diffusion models can generate realistic synthetic financial market data.
26 points by joegibbs | 7 comments
Hugging Face model card for TerrainSR, a diffusion model that upscales low-resolution terrain heightmaps (100m) to realistic higher-detail elevation data (10m).
13 points by gmays | 3 comments
arXiv paper introducing UniEvo-VL, an on-policy self-distillation framework where a multimodal model uses its own self-critiques to improve image generation.
49 points by schopra909 | 14 comments
Linum research blog post introducing Pyramid-JiT, a multiresolution prediction method to speed convergence of VAE-free pixel-space text-to-image models.
17 points by guanming0717 | 2 comments
GitHub repository for InstinctFlash, an AGPL-licensed high-performance serving runtime for robotics VLA and world-action models on RTX GPUs and Jetson Thor.
711 points by jmillikin | 189 comments
Qwen team blog post announcing Qwen-Image-2.1, a compact, efficient, unified image generation and editing model.
387 points by mattnewton | 43 comments
Krea's technical report introducing Krea 2, a family of open-weights image foundation models, covering data curation, architecture, training pipeline, and infrastructure.
365 points by janalsncm | 222 comments
GitHub repository for a Python CLI and library that removes visible (Gemini) and invisible (SynthID, C2PA, EXIF) AI watermarks from images.
418 points by dagenix | 115 comments
Project page for GenCAD, an MIT research model that generates CAD designs from images using transformer-based contrastive representations and diffusion priors.
378 points by meetpateltech | 160 comments
Qwen team blog post announcing Qwen-Image-2.0, an image generation model focused on professional infographics and photorealistic output.
421 points by antirez | 135 comments
GitHub repository by antirez for flux2.c, a pure C inference implementation of the FLUX.2 Klein image generation model.
368 points by doener | 145 comments
GitHub repository for Z-Image, a 6B-parameter image generation model from Alibaba's Tongyi-MAI team, released under Apache-2.0.
349 points by meetpateltech | 107 comments
Black Forest Labs blog post announcing FLUX.2, its new frontier image generation and editing model family, dated November 25, 2025.
406 points by lastdong | 150 comments
Project page for Microsoft's VibeVoice, an open-source TTS model generating up to 90 minutes of multi-speaker conversational audio using LLM and diffusion.
515 points by meetpateltech | 145 comments
Qwen Team blog post announcing Qwen-Image, an image generation foundation model notable for native text rendering, including Chinese text.
344 points by vmatsiiako | 98 comments
Krea blog post announcing the open-weights release of FLUX.1 Krea [dev], a guidance-distilled image model built with Black Forest Labs and compatible with FLUX.1-dev.
454 points by yousef_g | 73 comments
GitHub repository miniDiffusion, a from-scratch reimplementation of Stable Diffusion 3.5 in pure PyTorch, MIT licensed.
485 points by minimaxir | 113 comments
Black Forest Labs product page for FLUX.1 Kontext, a model suite for text-and-image-driven image generation and editing.
804 points by mdp2021 | 215 comments
Simon Willison's blog post on Google's Gemini Diffusion, a diffusion-based LLM announced at Google I/O, noting its speed of about 857 tokens per second.
362 points by __rito__ | 21 comments
Course page for MIT 6.S184, an introduction to flow matching and diffusion models built from stochastic differential equations, with notes, lectures, and labs.
455 points by francoisfleuret | 207 comments
Project page for DIAMOND, a NeurIPS 2024 spotlight paper on using diffusion models as world models for reinforcement learning agents in Atari games.
1087 points by jmorgan | 391 comments
Project page for GameNGen, a Google Research paper showing a diffusion model running a playable game in real time, with links to the paper, arXiv, and a demo video.
372 points by pouwerkerk | 85 comments
Project page for a UC Berkeley paper on applying diffusion over syntax trees so neural models iteratively edit programs guided by execution output.
710 points by ed | 123 comments
Stability AI announcement of Stable Video 3D (SV3D), a generative model based on Stable Video Diffusion for novel view synthesis and 3D generation from single images.
342 points by jxmorris12 | 36 comments
A tutorial blog post by Chenyang Yuan explaining diffusion models from scratch through a new theoretical perspective.
446 points by ed | 88 comments
Stability AI blog post announcing the research paper behind Stable Diffusion 3 and its Multimodal Diffusion Transformer (MMDiT) architecture.
900 points by reqo | 635 comments
Stability AI's announcement post for Stable Diffusion 3, its next-generation text-to-image generative model.
666 points by davidbarker | 154 comments
GitHub repository for Stability AI's Stable Cascade, a text-to-image diffusion model, with code for training and inference.
465 points by beefman | 199 comments
Demo page for Stability AI's Stable Audio model, showcasing generated long-form stereo music samples from text prompts, with links to the paper and code.
336 points by Flux159 | 61 comments
GitHub repository for StreamDiffusion, a pipeline-level solution for real-time interactive image generation with diffusion models.
758 points by beefman | 64 comments
Project page for Visual Anagrams, a University of Michigan research paper using pretrained diffusion models to generate multi-view optical illusions.
1279 points by roborovskis | 293 comments
Stability AI announcement of Stable Video Diffusion, its first open foundation model for generative video, released as a research preview with code and weights.
668 points by sandslides | 220 comments
GitHub repository for StyleTTS 2, an open-source text-to-speech model using style diffusion and adversarial training with large speech language models.
557 points by flavoredquark | 178 comments
Project page for RealFill, a Google Research and Cornell paper on reference-driven image completion using diffusion models.
426 points by varunshenoy | 113 comments
GitHub repository for Opendream, an extensible, MIT-licensed, layer-based web UI for Stable Diffusion image generation.
538 points by noob_eng | 101 comments
Overview page for fast.ai's Practical Deep Learning Part 2 course, a 30+ hour video course implementing Stable Diffusion from scratch.